mirror of https://github.com/hpcaitech/ColossalAI
Merge branch 'main' of https://github.com/hpcaitech/ColossalAI into main
jiaruifang
commit
cddb4b6f6f
|
|
@ -70,11 +70,6 @@
|
|||
<li><a href="#使用-Docker">使用 Docker</a></li>
|
||||
<li><a href="#社区">社区</a></li>
|
||||
<li><a href="#做出贡献">做出贡献</a></li>
|
||||
<li><a href="#快速预览">快速预览</a></li>
|
||||
<ul>
|
||||
<li><a href="#几行代码开启分布式训练">几行代码开启分布式训练</a></li>
|
||||
<li><a href="#构建一个简单的2维并行模型">构建一个简单的2维并行模型</a></li>
|
||||
</ul>
|
||||
<li><a href="#引用我们">引用我们</a></li>
|
||||
</ul>
|
||||
|
||||
|
|
@ -306,31 +301,6 @@ docker run -ti --gpus all --rm --ipc=host colossalai bash
|
|||
|
||||
<p align="right">(<a href="#top">返回顶端</a>)</p>
|
||||
|
||||
## 快速预览
|
||||
|
||||
### 几行代码开启分布式训练
|
||||
|
||||
```python
|
||||
parallel = dict(
|
||||
pipeline=2,
|
||||
tensor=dict(mode='2.5d', depth = 1, size=4)
|
||||
)
|
||||
```
|
||||
|
||||
### 几行代码开启异构训练
|
||||
|
||||
```python
|
||||
zero = dict(
|
||||
model_config=dict(
|
||||
tensor_placement_policy='auto',
|
||||
shard_strategy=TensorShardStrategy(),
|
||||
reuse_fp16_shard=True
|
||||
),
|
||||
optimizer_config=dict(initial_scale=2**5, gpu_margin_mem_ratio=0.2)
|
||||
)
|
||||
```
|
||||
|
||||
<p align="right">(<a href="#top">返回顶端</a>)</p>
|
||||
|
||||
## 引用我们
|
||||
|
||||
|
|
|
|||
31
README.md
31
README.md
|
|
@ -70,11 +70,6 @@
|
|||
<li><a href="#Use-Docker">Use Docker</a></li>
|
||||
<li><a href="#Community">Community</a></li>
|
||||
<li><a href="#contributing">Contributing</a></li>
|
||||
<li><a href="#Quick-View">Quick View</a></li>
|
||||
<ul>
|
||||
<li><a href="#Start-Distributed-Training-in-Lines">Start Distributed Training in Lines</a></li>
|
||||
<li><a href="#Write-a-Simple-2D-Parallel-Model">Write a Simple 2D Parallel Model</a></li>
|
||||
</ul>
|
||||
<li><a href="#Cite-Us">Cite Us</a></li>
|
||||
</ul>
|
||||
|
||||
|
|
@ -311,32 +306,6 @@ Thanks so much to all of our amazing contributors!
|
|||
|
||||
<p align="right">(<a href="#top">back to top</a>)</p>
|
||||
|
||||
## Quick View
|
||||
|
||||
### Start Distributed Training in Lines
|
||||
|
||||
```python
|
||||
parallel = dict(
|
||||
pipeline=2,
|
||||
tensor=dict(mode='2.5d', depth = 1, size=4)
|
||||
)
|
||||
```
|
||||
|
||||
### Start Heterogeneous Training in Lines
|
||||
|
||||
```python
|
||||
zero = dict(
|
||||
model_config=dict(
|
||||
tensor_placement_policy='auto',
|
||||
shard_strategy=TensorShardStrategy(),
|
||||
reuse_fp16_shard=True
|
||||
),
|
||||
optimizer_config=dict(initial_scale=2**5, gpu_margin_mem_ratio=0.2)
|
||||
)
|
||||
|
||||
```
|
||||
|
||||
<p align="right">(<a href="#top">back to top</a>)</p>
|
||||
|
||||
## Cite Us
|
||||
|
||||
|
|
|
|||
Binary file not shown.
|
Before Width: | Height: | Size: 3.8 MiB |
|
|
@ -1,21 +1,27 @@
|
|||
# ColoDiffusion
|
||||
*[ColoDiffusion](https://github.com/hpcaitech/ColoDiffusion) is a Faster Train implementation of the model [stable-diffusion](https://github.com/CompVis/stable-diffusion) from [Stability AI](https://stability.ai/)*
|
||||
# Stable Diffusion with Colossal-AI
|
||||
*[Colosssal-AI](https://github.com/hpcaitech/ColossalAI) provides a faster and lower cost solution for pretraining and
|
||||
fine-tuning for AIGC (AI-Generated Content) applications such as the model [stable-diffusion](https://github.com/CompVis/stable-diffusion) from [Stability AI](https://stability.ai/).*
|
||||
|
||||
We take advantage of Colosssal-AI to exploit multiple optimization strategies
|
||||
We take advantage of [Colosssal-AI](https://github.com/hpcaitech/ColossalAI) to exploit multiple optimization strategies
|
||||
, e.g. data parallelism, tensor parallelism, mixed precision & ZeRO, to scale the training to multiple GPUs.
|
||||
|
||||
|
||||
|
||||
|
||||
[Stable Diffusion](#stable-diffusion-v1) is a latent text-to-image diffusion
|
||||
## Stable Diffusion
|
||||
[Stable Diffusion](https://huggingface.co/CompVis/stable-diffusion) is a latent text-to-image diffusion
|
||||
model.
|
||||
Thanks to a generous compute donation from [Stability AI](https://stability.ai/) and support from [LAION](https://laion.ai/), we were able to train a Latent Diffusion Model on 512x512 images from a subset of the [LAION-5B](https://laion.ai/blog/laion-5b/) database.
|
||||
Similar to Google's [Imagen](https://arxiv.org/abs/2205.11487),
|
||||
this model uses a frozen CLIP ViT-L/14 text encoder to condition the model on text prompts.
|
||||
With its 860M UNet and 123M text encoder, the model is relatively lightweight and runs on a GPU with at least 10GB VRAM.
|
||||
See [this section](#stable-diffusion-v1) below and the [model card](https://huggingface.co/CompVis/stable-diffusion).
|
||||
|
||||
|
||||
<p id="diffusion_train" align="center">
|
||||
<img src="https://raw.githubusercontent.com/hpcaitech/public_assets/main/colossalai/img/diffusion_train.png" width=800/>
|
||||
</p>
|
||||
|
||||
[Stable Diffusion with Colossal-AI](https://github.com/hpcaitech/ColossalAI/tree/main/examples/images/diffusion) provides **6.5x faster training and pretraining cost saving, the hardware cost of fine-tuning can be almost 7X cheaper** (from RTX3090/4090 24GB to RTX3050/2070 8GB).
|
||||
|
||||
<p id="diffusion_demo" align="center">
|
||||
<img src="https://raw.githubusercontent.com/hpcaitech/public_assets/main/colossalai/img/diffusion_demo.png" width=800/>
|
||||
</p>
|
||||
|
||||
## Requirements
|
||||
A suitable [conda](https://conda.io/) environment named `ldm` can be created
|
||||
and activated with:
|
||||
|
|
@ -33,7 +39,7 @@ pip install transformers==4.19.2 diffusers invisible-watermark
|
|||
pip install -e .
|
||||
```
|
||||
|
||||
### Install ColossalAI
|
||||
### Install Colossal-AI
|
||||
|
||||
```
|
||||
git clone https://github.com/hpcaitech/ColossalAI.git
|
||||
|
|
@ -41,7 +47,7 @@ git checkout v0.1.10
|
|||
pip install .
|
||||
```
|
||||
|
||||
### Install colossalai lightning
|
||||
### Install Colossal-AI [Lightning](https://github.com/Lightning-AI/lightning)
|
||||
```
|
||||
git clone -b colossalai https://github.com/Fazziekey/lightning.git
|
||||
pip install .
|
||||
|
|
@ -74,16 +80,23 @@ you can change the trainging config in the yaml file
|
|||
## Comments
|
||||
|
||||
- Our codebase for the diffusion models builds heavily on [OpenAI's ADM codebase](https://github.com/openai/guided-diffusion)
|
||||
and [https://github.com/lucidrains/denoising-diffusion-pytorch](https://github.com/lucidrains/denoising-diffusion-pytorch).
|
||||
, [https://github.com/lucidrains/denoising-diffusion-pytorch](https://github.com/lucidrains/denoising-diffusion-pytorch),
|
||||
[Stable Diffusion](https://github.com/CompVis/stable-diffusion) and [Hugging Face](https://huggingface.co/CompVis/stable-diffusion).
|
||||
Thanks for open-sourcing!
|
||||
|
||||
- The implementation of the transformer encoder is from [x-transformers](https://github.com/lucidrains/x-transformers) by [lucidrains](https://github.com/lucidrains?tab=repositories).
|
||||
|
||||
- the implementation of [flash attention](https://github.com/HazyResearch/flash-attention) is from [HazyResearch](https://github.com/HazyResearch)
|
||||
- The implementation of [flash attention](https://github.com/HazyResearch/flash-attention) is from [HazyResearch](https://github.com/HazyResearch).
|
||||
|
||||
## BibTeX
|
||||
|
||||
```
|
||||
@article{bian2021colossal,
|
||||
title={Colossal-AI: A Unified Deep Learning System For Large-Scale Parallel Training},
|
||||
author={Bian, Zhengda and Liu, Hongxin and Wang, Boxiang and Huang, Haichen and Li, Yongbin and Wang, Chuanrui and Cui, Fan and You, Yang},
|
||||
journal={arXiv preprint arXiv:2110.14883},
|
||||
year={2021}
|
||||
}
|
||||
@misc{rombach2021highresolution,
|
||||
title={High-Resolution Image Synthesis with Latent Diffusion Models},
|
||||
author={Robin Rombach and Andreas Blattmann and Dominik Lorenz and Patrick Esser and Björn Ommer},
|
||||
|
|
|
|||
|
|
@ -1,144 +0,0 @@
|
|||
# Stable Diffusion v1 Model Card
|
||||
This model card focuses on the model associated with the Stable Diffusion model, available [here](https://github.com/CompVis/stable-diffusion).
|
||||
|
||||
## Model Details
|
||||
- **Developed by:** Robin Rombach, Patrick Esser
|
||||
- **Model type:** Diffusion-based text-to-image generation model
|
||||
- **Language(s):** English
|
||||
- **License:** [Proprietary](LICENSE)
|
||||
- **Model Description:** This is a model that can be used to generate and modify images based on text prompts. It is a [Latent Diffusion Model](https://arxiv.org/abs/2112.10752) that uses a fixed, pretrained text encoder ([CLIP ViT-L/14](https://arxiv.org/abs/2103.00020)) as suggested in the [Imagen paper](https://arxiv.org/abs/2205.11487).
|
||||
- **Resources for more information:** [GitHub Repository](https://github.com/CompVis/stable-diffusion), [Paper](https://arxiv.org/abs/2112.10752).
|
||||
- **Cite as:**
|
||||
|
||||
@InProceedings{Rombach_2022_CVPR,
|
||||
author = {Rombach, Robin and Blattmann, Andreas and Lorenz, Dominik and Esser, Patrick and Ommer, Bj\"orn},
|
||||
title = {High-Resolution Image Synthesis With Latent Diffusion Models},
|
||||
booktitle = {Proceedings of the IEEE/CVF Conference on Computer Vision and Pattern Recognition (CVPR)},
|
||||
month = {June},
|
||||
year = {2022},
|
||||
pages = {10684-10695}
|
||||
}
|
||||
|
||||
# Uses
|
||||
|
||||
## Direct Use
|
||||
The model is intended for research purposes only. Possible research areas and
|
||||
tasks include
|
||||
|
||||
- Safe deployment of models which have the potential to generate harmful content.
|
||||
- Probing and understanding the limitations and biases of generative models.
|
||||
- Generation of artworks and use in design and other artistic processes.
|
||||
- Applications in educational or creative tools.
|
||||
- Research on generative models.
|
||||
|
||||
Excluded uses are described below.
|
||||
|
||||
### Misuse, Malicious Use, and Out-of-Scope Use
|
||||
_Note: This section is taken from the [DALLE-MINI model card](https://huggingface.co/dalle-mini/dalle-mini), but applies in the same way to Stable Diffusion v1_.
|
||||
|
||||
The model should not be used to intentionally create or disseminate images that create hostile or alienating environments for people. This includes generating images that people would foreseeably find disturbing, distressing, or offensive; or content that propagates historical or current stereotypes.
|
||||
|
||||
#### Out-of-Scope Use
|
||||
The model was not trained to be factual or true representations of people or events, and therefore using the model to generate such content is out-of-scope for the abilities of this model.
|
||||
|
||||
#### Misuse and Malicious Use
|
||||
Using the model to generate content that is cruel to individuals is a misuse of this model. This includes, but is not limited to:
|
||||
|
||||
- Generating demeaning, dehumanizing, or otherwise harmful representations of people or their environments, cultures, religions, etc.
|
||||
- Intentionally promoting or propagating discriminatory content or harmful stereotypes.
|
||||
- Impersonating individuals without their consent.
|
||||
- Sexual content without consent of the people who might see it.
|
||||
- Mis- and disinformation
|
||||
- Representations of egregious violence and gore
|
||||
- Sharing of copyrighted or licensed material in violation of its terms of use.
|
||||
- Sharing content that is an alteration of copyrighted or licensed material in violation of its terms of use.
|
||||
|
||||
## Limitations and Bias
|
||||
|
||||
### Limitations
|
||||
|
||||
- The model does not achieve perfect photorealism
|
||||
- The model cannot render legible text
|
||||
- The model does not perform well on more difficult tasks which involve compositionality, such as rendering an image corresponding to “A red cube on top of a blue sphere”
|
||||
- Faces and people in general may not be generated properly.
|
||||
- The model was trained mainly with English captions and will not work as well in other languages.
|
||||
- The autoencoding part of the model is lossy
|
||||
- The model was trained on a large-scale dataset
|
||||
[LAION-5B](https://laion.ai/blog/laion-5b/) which contains adult material
|
||||
and is not fit for product use without additional safety mechanisms and
|
||||
considerations.
|
||||
- No additional measures were used to deduplicate the dataset. As a result, we observe some degree of memorization for images that are duplicated in the training data.
|
||||
The training data can be searched at [https://rom1504.github.io/clip-retrieval/](https://rom1504.github.io/clip-retrieval/) to possibly assist in the detection of memorized images.
|
||||
|
||||
### Bias
|
||||
While the capabilities of image generation models are impressive, they can also reinforce or exacerbate social biases.
|
||||
Stable Diffusion v1 was primarily trained on subsets of [LAION-2B(en)](https://laion.ai/blog/laion-5b/),
|
||||
which consists of images that are limited to English descriptions.
|
||||
Texts and images from communities and cultures that use other languages are likely to be insufficiently accounted for.
|
||||
This affects the overall output of the model, as white and western cultures are often set as the default. Further, the
|
||||
ability of the model to generate content with non-English prompts is significantly worse than with English-language prompts.
|
||||
Stable Diffusion v1 mirrors and exacerbates biases to such a degree that viewer discretion must be advised irrespective of the input or its intent.
|
||||
|
||||
|
||||
## Training
|
||||
|
||||
**Training Data**
|
||||
The model developers used the following dataset for training the model:
|
||||
|
||||
- LAION-5B and subsets thereof (see next section)
|
||||
|
||||
**Training Procedure**
|
||||
Stable Diffusion v1 is a latent diffusion model which combines an autoencoder with a diffusion model that is trained in the latent space of the autoencoder. During training,
|
||||
|
||||
- Images are encoded through an encoder, which turns images into latent representations. The autoencoder uses a relative downsampling factor of 8 and maps images of shape H x W x 3 to latents of shape H/f x W/f x 4
|
||||
- Text prompts are encoded through a ViT-L/14 text-encoder.
|
||||
- The non-pooled output of the text encoder is fed into the UNet backbone of the latent diffusion model via cross-attention.
|
||||
- The loss is a reconstruction objective between the noise that was added to the latent and the prediction made by the UNet.
|
||||
|
||||
We currently provide the following checkpoints:
|
||||
|
||||
- `sd-v1-1.ckpt`: 237k steps at resolution `256x256` on [laion2B-en](https://huggingface.co/datasets/laion/laion2B-en).
|
||||
194k steps at resolution `512x512` on [laion-high-resolution](https://huggingface.co/datasets/laion/laion-high-resolution) (170M examples from LAION-5B with resolution `>= 1024x1024`).
|
||||
- `sd-v1-2.ckpt`: Resumed from `sd-v1-1.ckpt`.
|
||||
515k steps at resolution `512x512` on [laion-aesthetics v2 5+](https://laion.ai/blog/laion-aesthetics/) (a subset of laion2B-en with estimated aesthetics score `> 5.0`, and additionally
|
||||
filtered to images with an original size `>= 512x512`, and an estimated watermark probability `< 0.5`. The watermark estimate is from the [LAION-5B](https://laion.ai/blog/laion-5b/) metadata, the aesthetics score is estimated using the [LAION-Aesthetics Predictor V2](https://github.com/christophschuhmann/improved-aesthetic-predictor)).
|
||||
- `sd-v1-3.ckpt`: Resumed from `sd-v1-2.ckpt`. 195k steps at resolution `512x512` on "laion-aesthetics v2 5+" and 10\% dropping of the text-conditioning to improve [classifier-free guidance sampling](https://arxiv.org/abs/2207.12598).
|
||||
- `sd-v1-4.ckpt`: Resumed from `sd-v1-2.ckpt`. 225k steps at resolution `512x512` on "laion-aesthetics v2 5+" and 10\% dropping of the text-conditioning to improve [classifier-free guidance sampling](https://arxiv.org/abs/2207.12598).
|
||||
|
||||
- **Hardware:** 32 x 8 x A100 GPUs
|
||||
- **Optimizer:** AdamW
|
||||
- **Gradient Accumulations**: 2
|
||||
- **Batch:** 32 x 8 x 2 x 4 = 2048
|
||||
- **Learning rate:** warmup to 0.0001 for 10,000 steps and then kept constant
|
||||
|
||||
## Evaluation Results
|
||||
Evaluations with different classifier-free guidance scales (1.5, 2.0, 3.0, 4.0,
|
||||
5.0, 6.0, 7.0, 8.0) and 50 PLMS sampling
|
||||
steps show the relative improvements of the checkpoints:
|
||||
|
||||
|
||||
|
||||
Evaluated using 50 PLMS steps and 10000 random prompts from the COCO2017 validation set, evaluated at 512x512 resolution. Not optimized for FID scores.
|
||||
|
||||
## Environmental Impact
|
||||
|
||||
**Stable Diffusion v1** **Estimated Emissions**
|
||||
Based on that information, we estimate the following CO2 emissions using the [Machine Learning Impact calculator](https://mlco2.github.io/impact#compute) presented in [Lacoste et al. (2019)](https://arxiv.org/abs/1910.09700). The hardware, runtime, cloud provider, and compute region were utilized to estimate the carbon impact.
|
||||
|
||||
- **Hardware Type:** A100 PCIe 40GB
|
||||
- **Hours used:** 150000
|
||||
- **Cloud Provider:** AWS
|
||||
- **Compute Region:** US-east
|
||||
- **Carbon Emitted (Power consumption x Time x Carbon produced based on location of power grid):** 11250 kg CO2 eq.
|
||||
|
||||
## Citation
|
||||
@InProceedings{Rombach_2022_CVPR,
|
||||
author = {Rombach, Robin and Blattmann, Andreas and Lorenz, Dominik and Esser, Patrick and Ommer, Bj\"orn},
|
||||
title = {High-Resolution Image Synthesis With Latent Diffusion Models},
|
||||
booktitle = {Proceedings of the IEEE/CVF Conference on Computer Vision and Pattern Recognition (CVPR)},
|
||||
month = {June},
|
||||
year = {2022},
|
||||
pages = {10684-10695}
|
||||
}
|
||||
|
||||
*This model card was written by: Robin Rombach and Patrick Esser and is based on the [DALL-E Mini model card](https://huggingface.co/dalle-mini/dalle-mini).*
|
||||
Loading…
Reference in New Issue